This repository provides the simplest tutorial code for AIGC researchers to use Lora in just a few lines. Using this handbook, you can easily play with any Lora model from active communities such as Huggingface and cititai.
Now, we also support ControlNet-for-Diffusers, T2I-Adapter-for-Diffusers.
Low-Rank Adaptation of Large Language Models (LoRA) is developed by Microsoft to reduce the number of trainable parameters by learning pairs of rank-decompostion matrices while freezing the original weights. Lora attemptes to fine-tune the "residual" of the model instead of the entire model: i.e., train the
Where
This training trick is quite useful for fune-tuning customized models on a large general base model. Various text to image models have been developed built on the top of the official Stable Diffusion. Now, with Lora, you can efficiently train your own model with much less resources.
Safetensors is a new simple format for storing tensors safely (as opposed to pickle) released by Hugging Face and that is still fast (zero-copy). For its efficiency, many stable diffusion models, especially Lora models are released in safetensors format. You can find more its advantages from huggingface/safetensors and install it via pip install.
pip install safetensorsIn this tutorial, we show to load or insert pre-trained Lora into diffusers framework. Many interesting projects can be found in Huggingface and cititai, but mostly in stable-diffusion-webui framework, which is not convenient for advanced developers. We highly motivated by cloneofsimo/lora about loading, merging, and interpolating trained LORAs. We mainly discuss models in safetensors format which is not well compatible with diffusers.
A full model includes all modules needed (base model with or without Lora layers), they are usually stored in .ckpt or .safetensors format. We provide two examples below to show you how to use on hand.
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stabilityai/stable-diffusion-2-1 from Huggingface.
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dreamshaper from Civitai.
You can download .ckpt or .safetensors file only. Although diffusers does not support loading them directly, they do provide the converting script. First download diffusers to local.
git clone https://github.com/huggingface/diffuserscd ./diffusers
# assume you have downloaded xxx.safetensors, it will out save_dir in diffusers format.
python ./scripts/convert_original_stable_diffusion_to_diffusers.py --checkpoint_path xxx.safetensors --dump_path save_dir --from_safetensors
# assume you have downloaded xxx.ckpt, it will out save_dir in diffusers format.
python ./scripts/convert_original_stable_diffusion_to_diffusers.py --checkpoint_path xxx.ckpt --dump_path save_dirThen, you can load the model
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained(save_dir,torch_dtype=torch.float32)For now, diffusers cannot support load weights in Lora (usually in .safetensor format) . Here we show our attempts in an inelegant style. We also provide one example.
Note that the size of file is much smaller than full model, as it only contains extra Lora weights. In the case, we have to load the base model. It is also fine to just load stable-diffusion 1.5 as base, but to get satisfied results, it is recommanded to download suggested base model.
Our method is very straightforward: take out weight from .safetensor, and merge lora weight into a diffusers supported weight. We don't convert .safetensor into other format, we update the weight of base model instead.
Our script should work fine with most of models from Huggingface and cititai, if not, you can also modify the code on your own. Believe me, it is really simple and you can make it.
# the default mergering ratio is 0.75, you can manually set it
python convert_lora_safetensor_to_diffusers.pyWe have made a PR for diffusers on this issue, where we further warp the convering function so that it is more flexible. You can directly check it if you cannot wait. It shall be merged into diffusers soon!
Diffusers has provide a simple train_text_to_image_lora.py to train your on Lora model. Please follow its instruction to install requirements.
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch --mixed_precision="fp16" train_text_to_image_lora.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--dataset_name=$DATASET_NAME --caption_column="text" \
--resolution=512 --random_flip \
--train_batch_size=1 \
--num_train_epochs=100 --checkpointing_steps=5000 \
--learning_rate=1e-04 --lr_scheduler="constant" --lr_warmup_steps=0 \
--seed=42 \
--output_dir="sd-pokemon-model-lora" \
--validation_prompt="cute dragon creature" --report_to="wandb"Once you have trained a model using above command, the inference can be done simply using the StableDiffusionPipeline after loading the trained LoRA weights. You need to pass the output_dir for loading the LoRA weights which, in this case, is sd-pokemon-model-lora.
import torch
from diffusers import StableDiffusionPipeline
model_path = "your_path/sd-model-finetuned-lora-t4"
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipe.unet.load_attn_procs(model_path)
pipe.to("cuda")
prompt = "A pokemon with green eyes and red legs."
image = pipe(prompt, num_inference_steps=30, guidance_scale=7.5).images[0]
image.save("pokemon.png")For now, diffusers only supports train LoRA for UNet. We have supported and made a PR, if you need it, please check with our PR or open an issue.
ColossalAI supports LoRA already. We only need modify a few lines on the top of train_dreambooth_colossalai.py. This example is for dreambooth, but you can easily adopt it regular text to image training. The generated LoRA weights are only for attention layers in UNet. If you want to support text encoder too, please use acceletate framework in diffusers, as ColossalAI does not support multiple models yet.
from diffusers.loaders import AttnProcsLayers
from diffusers.models.cross_attention import LoRACrossAttnProcessor
# attention here! It is necessaray to init unet under ColoInitContext, not just lora layers
with ColoInitContext(device=get_current_device()):
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path,
subfolder="unet",
revision=args.revision,
low_cpu_mem_usage=False
)
unet.requires_grad_(False)
# Set correct lora layers
lora_attn_procs = {}
for name in unet.attn_processors.keys():
cross_attention_dim = None if name.endswith("attn1.processor") else unet.config.cross_attention_dim
if name.startswith("mid_block"):
hidden_size = unet.config.block_out_channels[-1]
elif name.startswith("up_blocks"):
block_id = int(name[len("up_blocks.")])
hidden_size = list(reversed(unet.config.block_out_channels))[block_id]
elif name.startswith("down_blocks"):
block_id = int(name[len("down_blocks.")])
hidden_size = unet.config.block_out_channels[block_id]
lora_attn_procs[name] = LoRACrossAttnProcessor(
hidden_size=hidden_size, cross_attention_dim=cross_attention_dim
)
unet.set_attn_processor(lora_attn_procs)
lora_layers = AttnProcsLayers(unet.attn_processors)
# DDP
unet = gemini_zero_dpp(unet, args.placement)
# config optimizer for colossalai zero, set initial_scale to large value to avoid underflow
optimizer = GeminiAdamOptimizer(unet,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
initial_scale=2**16,
clipping_norm=args.max_grad_norm)Here we go, the only thing is the initialization way of UNet. To save LoRA weights only,
torch_unet = get_static_torch_model(unet)
if gpc.get_local_rank(ParallelMode.DATA) == 0:
torch_unet = torch_unet.to(torch.float32)
torch_unet.save_attn_procs(save_path)Then, do inference
from diffusers import StableDiffusionPipeline
import torch
model_path = "sd-model-finetuned-lora"
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipe.unet.load_attn_procs(model_path)
pipe.to("cuda")
prompt = "A pokemon with green eyes and red legs."
image = pipe(prompt, num_inference_steps=30, guidance_scale=7.5).images[0]
image.save("pokemon.png")You may find that the generated LoRA weight is only about 3MB size, this is because of the default setting. To increase the size, you can manually set the rank (dimension for low rank decomposition) for LoRA layers.
lora_attn_procs[name] = LoRACrossAttnProcessor(hidden_size=hidden_size, cross_attention_dim=cross_attention_dim, rank=128)Then, the LoRA weights will be about 100-200MB size. Be aware that LoRA layers are easy to overfit, generally speaking, it should be enough to train only 100 - 2000 steps on small datasets (less than 1K images) with batch size = 64.
(1) Can I manually adjust the weight of LoRA when merging?
Yes, the alpha here is the weight for LoRA. We have submitted a PR to diffusers where we provide the flexible warpped function.
(2) Can I only convert LoRA (.safetensors) into other formats that diffusers supported?
You can but we don't suggest, see this issues. There are many limitations. For example, our script cannot generalize to all .safetensors because some of them have different naming. Besides, current diffusers framework only supports adding LoRA into UNet's attention layers, while many .safetensors from civitai contain LoRA weights for other modules like text encoder. But LoRA for text encoder should be supported soon.
(3) Can I mix more than one LoRA model?
Yes, the only thing is to merge twice. But please carefully set the alpha (weight of LoRA), model degrades if alpha is too large.
(4) What's the motivation of this project?
We find there are many incredible models in civitai platform, but most of LoRA weights are in safetensors format, which is not convenient for diffusers users. Thus, we write a converting script so that you can use these LoRAs in diffusers. Be aware that we are not target for stable-diffusion-webui, which is already very mature but has totally different API as diffusers.